PLEASE HELP


A frustum of a regular square pyramid has bases with sides of lengths 6 and 10. The height of the frustum is 12.


Find the volume of the frustum?


Find the surface area of the frustum?

The pizza from Jaco, Costa Rica, with a 27.8-centimeter diameter, has approximately 603.42 square centimeters of pizza.

To calculate the number of square centimeters of pizza, we need to determine the area of the circle using the formula A = πr^2, where A is the area and r is the radius of the circle.

Finding the radius:

The diameter of the pizza from Jaco, Costa Rica, is given as 27.8 centimeters. To find the radius, we divide the diameter by 2:

Radius = Diameter / 2 = 27.8 cm / 2 = 13.9 cm

Calculating the area:

Now that we have the radius, we can substitute it into the formula:

A = πr^2 = π * (13.9 cm)^2

Using the value of π (pi) as approximately 3.14159, we can calculate the area:

A ≈ 3.14159 * (13.9 cm)^2 ≈ 3.14159 * 192.21 cm^2 ≈ 603.42 cm^2

Therefore, the pizza from Jaco, Costa Rica, with a 27.8-centimeter diameter, has approximately 603.42 square centimeters of pizza.

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Answer:Answer:The equation that could be used to determine the length of each of the four pieces of wood he cut is 3.5 = 0.6 + 4x and the solution is x = 0.725

Part A: Write an equation that could be used to determine the length of each of the four pieces of wood he cut.

Represent the length of the four pieces with x

So, the given parameters are:

Initial length = 3.5 feet

Remaining length = 0.6 feet

Number of pieces = 4

The equation that could be used to determine the length of each of the four pieces of wood he cut is represented as:

Initial length = Remaining length + Number of pieces * x

This gives

3.5 = 0.6 + 4x

Hence, the equation that could be used to determine the length of each of the four pieces of wood he cut is 3.5 = 0.6 + 4x

Part B: Explain how you know the equation from Part A is correct.

The equation in part (A) is correct because it can be used to determine the length of each of the four pieces of wood he cut

Part C: Solve the equation from Part A.

In part A, we have:

3.5 = 0.6 + 4x

Subtract 0.6 from both sides

2.9 = 4x

Divide both sides by 4

x = 0.725

Hence, the solution is x = 0.725

Step-by-step explanation:

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The speed of the particle along the path c(t) = (cos(4t), sin(4t), 3t) at time t = 0 is E. 5.

To find the speed of the particle at time t = 0, we need to find the magnitude of its velocity vector at that time. The speed at which an object's position changes is represented by a velocity vector. A velocity vector's magnitude indicates an object's speed, whereas the vector's direction indicates its direction. According to the vector addition tenets, velocity vectors can be added or deleted.
The velocity vector is given by the derivative of the position vector:
c'(t) = (-4sin(4t), 4cos(4t), 3)

At t = 0, we have:
c'(0) = (-4sin(0), 4cos(0), 3) = (0, 4, 3)

The magnitude of this vector is:
|c'(0)| = sqrt(0^2 + 4^2 + 3^2) = sqrt(25) = 5

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The set s consisting of pairwise orthogonal non-zero vectors in Rn is linearly independent.

To prove that the set s is linearly independent, we need to show that the only solution to the equation:

c1v1 + c2v2 + ... + cnvn = 0

is the trivial solution c1 = c2 = ... = cn = 0.

Suppose there exists a non-trivial solution to the above equation, i.e., there exists some non-zero vector c = (c1, c2, ..., cn) such that:

c1v1 + c2v2 + ... + cnvn = 0

Then, taking the dot product of both sides with vi, we get:

(ci vi)· vi = 0

since the dot product of any two orthogonal vectors is zero.

Thus, we have:

civi · vi = 0

or

civi² = 0

since vi·vi = ||vi||² ≠ 0, as each vector is nonzero.

Since each vector in s is nonzero, this implies that ci = 0 for all i, since the square of any nonzero scalar is nonzero. Therefore, the only solution to the equation c1v1 + c2v2 + ... + cnvn = 0 is the trivial solution c1 = c2 = ... = cn = 0.

Thus, the set s is linearly independent

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The problem of determining the satisfiability of boolean formula in disjunctive normal form (DNF) is known as the DNF-SAT problem. This problem can be solved in polynomial time using an algorithm called the resolution algorithm. The resolution algorithm works by repeatedly applying the resolution rule to simplify the formula until it is either determined to be satisfiable or unsatisfiable.

DNF is a standard form of representing boolean formulas, where the formula is expressed as a disjunction of conjunctions of literals. The DNF-SAT problem involves determining whether there exists an assignment of truth values to the variables in the formula that makes the formula true.

The resolution algorithm is a complete and sound method for solving the DNF-SAT problem. It works by iteratively applying the resolution rule, which allows two clauses to be combined into a new clause that is a logical consequence of the original clauses. The algorithm continues until either a contradiction is reached (meaning the formula is unsatisfiable) or until the formula is simplified to a single clause (meaning the formula is satisfiable).

In conclusion, the DNF-SAT problem is polynomial-time solvable using the resolution algorithm. This is an important result in computational complexity theory because it shows that some boolean formula problems can be solved efficiently, which has implications for the development of algorithms in other fields, such as artificial intelligence and optimization.

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Answer:

see explanation

Step-by-step explanation:

180° on the line is a straight angle

The question is given as: If a test of the model shows that it holds 8,000 ounces, will the bridge hold 1 ton? 8,000 ounces on the model is equal to _ ounces on the actual bridge. Convert ounces to pounds. The actual bridge can hold _ pounds. Therefore, the bridge will/will not hold 1 ton.

In order to answer the question, let's first convert the 8,000 ounces to pounds as follows: 1 pound = 16 ounces. Therefore, 1 ounce = 1/16 pounds.

Now, 8,000 ounces = 8,000/16 = 500 pounds8,000 ounces on the model is equal to 500 pounds on the actual bridge.

Now, let's find out how many pounds one ton is: 1 ton = 2,000 pounds.

Therefore, the actual bridge can hold 2,000 pounds.

Thus, since 2,000 pounds is greater than 500 pounds, the bridge will hold 1 ton.

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To evaluate the line integral along the path C given by x = 2t, y = 4t, where 0 ≤ t ≤ 1, we can follow these steps:

1. Rewrite the given integral in terms of t using the parameterization of the path: C: x = 2t, y = 4t.
2. Compute the derivatives dx/dt and dy/dt.
3. Substitute the parameterization and derivatives into the line integral.
4. Evaluate the integral over the specified interval.

Step 1:
The integral in terms of t is: ∫(3y² dy)

Step 2:
dx/dt = 2
dy/dt = 4

Step 3:
Substitute the parameterization and derivatives:
∫(3(4t)² * 4 dt) over the interval [0, 1]

Step 4:
Evaluate the integral:
∫(3 * 16t² * 4 dt) from 0 to 1
= 192 ∫(t² dt) from 0 to 1

Now, integrate and evaluate the integral:
= 192 * [1/3 * t^3] from 0 to 1
= 192 * (1/3 * 1^3 - 1/3 * 0^3)
= 64

So, the value of the line integral along the path C is 64.

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Two medical technologies that have helped to combat the spread of diseases in cities include:

Artificial intelligence

Telemedicine

There are different forms of medical technology that have helped in combatting diseases in cities. Some of these include artificial intelligence and telemedicine. Artificial intelligence has helped to combat diseases because the medical records of patients can be easily tracked and used in suggesting diagnoses to medical doctors.

Telemedicine has also helped as technological devices are used to deliver healthcare services in a fast and efficient manner.

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Given that a soccer team has 12 players. It is known that only 6 players are allowed on the field at a time. How many different groups of players can be on the field at a time?To determine the number of different groups of players that can be on the field at a time, we need to apply combination formula because the order does not matter when choosing the 6 players from the total of 12 players.

The formula for combination is given by:[tex]C(n, r) = \frac{n!}{r!(n - r)!}[/tex] where C is the number of combinations possible, n is the total number of items, and r is the number of items being chosen.Using the combination formula to calculate the number of different groups of players that can be on the field at a time[tex]C(12, 6) = \frac{12!}{6!(12 - 6)!}$$$$C(12, 6) = \frac{12!}{6!6!}$$$$C(12, 6) = \frac{12 × 11 × 10 × 9 × 8 × 7}{6 × 5 × 4 × 3 × 2 × 1 × 6 × 5 × 4 × 3 × 2 × 1}$$$$C(12, 6) = 924[/tex]

Therefore, there are 924 different groups of players that can be on the field at a time.

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The predicted Durante humor score for an individual with an IQ score of 110 is 178.

Based on the information given, we know that there is a strong, positive, linear relationship between an individual's IQ score and their sense of humor. Additionally, the psychologist has found a correlation coefficient of 0.81 between the two variables.
To predict the Durante humor score of an individual with an IQ score of 110, we can use the formula for a simple linear regression:
y = b0 + b1x
where y is the predicted Durante humor score, x is the IQ score, b0 is the intercept, and b1 is the slope of the regression line.
To find the intercept and slope, we need to use the sample means and standard deviations provided:
b1 = r * (Sy / Sx)
where r is the correlation coefficient and Sy and Sx are the standard deviations of the Durante humor scores and IQ scores, respectively.
b0 = ybar - b1 * xbar
where ybar and xbar are the sample means of the Durante humor scores and IQ scores, respectively.
Plugging in the values, we get:
b1 = 0.81 * (24 / 12) = 1.62
b0 = 140 - 1.62 * 105 = -3.1
Now we can use these values to predict the Durante humor score of an individual with an IQ score of 110:
y = -3.1 + 1.62 * 110 = 177.9
Therefore, the predicted Durante humor score for an individual with an IQ score of 110 is 178.

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(a) Let's analyze the relation R defined as XRy if y can be obtained from x by swapping any two bits.

To determine if R is an equivalence relation, we need to check three conditions: reflexivity, symmetry, and transitivity.

Reflexivity: For any x in D, we need to check if xRx holds true.

In this case, swapping any two bits of x with itself will result in the same value x. Therefore, xRx holds true for all x in D.

Symmetry: For any x and y in D, if xRy holds true, then yRx should also hold true.

Swapping any two bits of x to obtain y and then swapping the same two bits of y will result in x again. Thus, if xRy is true, yRx is also true.

Transitivity: For any x, y, and z in D, if xRy and yRz hold true, then xRz should also hold true.

If we can obtain y from x by swapping two bits and obtain z from y by swapping two bits, we can perform both swaps together to obtain z from x. Therefore, if xRy and yRz are true, xRz is also true.

Since the relation R satisfies all three conditions (reflexivity, symmetry, and transitivity), we can conclude that R is an equivalence relation.

(b) Let's analyze the relation R defined as XRy if y can be obtained from x by reordering the bits in any way.

To determine if R is an equivalence relation, we again need to check the three conditions: reflexivity, symmetry, and transitivity.

Reflexivity: For any x in D, we need to check if xRx holds true.

Reordering the bits of x in any way will still result in x itself. Therefore, xRx holds true for all x in D.

Symmetry: For any x and y in D, if xRy holds true, then yRx should also hold true.

Reordering the bits of x to obtain y and then reordering the bits of y will still result in x. Thus, if xRy is true, yRx is also true.

Transitivity: For any x, y, and z in D, if xRy and yRz hold true, then xRz should also hold true.

If we can obtain y from x by reordering the bits and obtain z from y by reordering the bits, we can combine the two reorderings to obtain z from x. Therefore, if xRy and yRz are true, xRz is also true.

Since the relation R satisfies all three conditions (reflexivity, symmetry, and transitivity), we can conclude that R is an equivalence relation.

In summary:

Relation R in part (a) is an equivalence relation.

Relation R in part (b) is also an equivalence relation.

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The line integral ∫cyds is equal to 7 + (2/3).

To evaluate the line integral ∫cyds, where the curve C is defined by the line segment from point (0, 1, 1) to point (2, 2, 3), and the vector field F(x, y) = (2x + 3)i + (3x + 2y)j, we need to parameterize the curve and calculate the dot product of the vector field and the tangent vector.

Let's start by finding the parameterization of the line segment C.

The equation of the line passing through the two points can be written as:

x = 2t

y = 1 + t

z = 1 + 2t

where t ranges from 0 to 1.

The tangent vector to the curve C can be found by differentiating the parameterization with respect to t:

r'(t) = (2, 1, 2)

Now, let's calculate the line integral using the parameterization of the curve and the vector field:

∫cyds = ∫(0 to 1) F(x, y) ⋅ r'(t) dt

Substituting the values for F(x, y) and r'(t), we have:

∫cyds = ∫(0 to 1) [(2(2t) + 3)(2) + (3(2t) + 2(1 + t))(1)] dt

Simplifying further, we get:

∫cyds = ∫(0 to 1) (4t + 3 + 6t + 2 + 2t + 2t^2) dt

∫cyds = ∫(0 to 1) (10t + 2 + 2t^2) dt

Integrating term by term, we have:

∫cyds = [5t^2 + 2t^3 + (2/3)t^3] evaluated from 0 to 1

Evaluating the integral, we get:

∫cyds = [5(1)^2 + 2(1)^3 + (2/3)(1)^3] - [5(0)^2 + 2(0)^3 + (2/3)(0)^3]

∫cyds = 5 + 2 + (2/3) - 0 - 0 - 0

∫cyds = 7 + (2/3)

Therefore, the line integral ∫cyds is equal to 7 + (2/3).

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Based on the given information that AD = DC = AC = a, we can conclude that the triangle is a 30-60-90 right triangle.

To prove that the triangle is a 30-60-90 right triangle, we can use the properties of this specific triangle.

In a 30-60-90 triangle, the sides are in a specific ratio. Let's denote the length of the shortest side as "a". Then the other sides can be determined as follows:

The length of the side opposite the 30-degree angle is "a".

The length of the side opposite the 60-degree angle is "a√3".

The length of the hypotenuse (the longest side) is "2a".

Given that AD = DC = AC = a, we can conclude that the triangle is indeed a 30-60-90 right triangle.

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--The given question is incomplete, the complete question is given below " Prove 30 - 60 - 90 Right Triangle

using the given data from the figure, the indicated length of the sides of the triangle

AD= DC= AC= a"--

Answer:

62

Step-by-step explanation:

7/13 as a decimal to the hundredths place is 0.54 and the remainder as a fraction is 7/13.

7/13 as a decimal to the hundredths place and the remainder as a fraction

In order to convert 7/13 to a decimal, we will divide 7 by 13.

Using long division, we get7 ÷ 13 = 0.53846153846...To the nearest hundredth, we round up to 0.54.

Hence, 7/13 as a decimal to the hundredths place is 0.54.

To find the remainder as a fraction, we subtract the product of the quotient and divisor from the dividend. Then, we simplify the fraction as much as possible.

Remainder = Dividend - Quotient x DivisorRemainder = 7 - 0 x 13

Remainder = 7/13

Therefore, 7/13 as a decimal to the hundredths place is 0.54 and the remainder as a fraction is 7/13.

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The values of x that make the function undefined are x = 2 and x = -2. These values make the denominator equal to zero, causing the function to be undefined. Note that x = 3 does not make the denominator zero, so the function is defined for x = 3.

To determine the values that would make the function f(x) = (3-x) / (x^2 - 4) undefined, we need to find the values of x for which the denominator of the fraction becomes zero.

The denominator of the function is x^2 - 4. To find the values of x that make the denominator zero, we'll set x^2 - 4 equal to zero and solve for x:

x^2 - 4 = 0

We can factor this expression as a difference of squares:

(x + 2)(x - 2) = 0

Now, we'll solve for x by setting each factor equal to zero:

1) x + 2 = 0
  x = -2

2) x - 2 = 0
  x = 2

So, the values of x that make the function undefined are x = 2 and x = -2. These values make the denominator equal to zero, causing the function to be undefined. Note that x = 3 does not make the denominator zero, so the function is defined for x = 3.

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To construct a value-weighted portfolio, we need to allocate funds based on the market capitalization of each stock. The total market cap of the four stocks is $159 million. Therefore, OGG represents 32.7%, HNL represents 47.8%, KOA represents 11.9%, and LIH represents 7.5% of the total market cap. If we have $501,000 to invest, we should invest $163,710 in OGG, $239,430 in HNL, $59,490 in KOA, and $37,370 in LIH.

A value-weighted portfolio is a strategy that allocates funds based on the market capitalization of each stock. It means investing more in companies with a higher market capitalization and less in companies with a lower market capitalization. In this case, we calculate the percentage of each stock's market capitalization to the total market capitalization of all four stocks and allocate funds accordingly.

To construct a value-weighted portfolio of the four stocks, we should allocate funds based on the market capitalization of each stock. In this case, we allocate funds in the proportion of 32.7%, 47.8%, 11.9%, and 7.5% for OGG, HNL, KOA, and LIH, respectively. This ensures that we invest more in companies with a higher market capitalization and less in companies with a lower market capitalization.

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Answer:

x=17 degrees

Step-by-step explanation:

All 3 angles = 180 degrees

So 90 + 54 + (x+19) = 180

Combine like terms

163 + x = 180

Subtract 163 from both sides

x = 180-163

x = 17

A linearly independent system of vectors v1, v2, ..., vn in a vector space v is a basis if and only if the number of vectors, n, is equal to the dimension of v.

To prove that a linearly independent system of vectors v1, v2, ..., vn in a vector space v is a basis if and only if n = dim v, we need to show both directions of the statement.

If the system of vectors is a basis, then n = dim v:

Suppose the system of vectors v1, v2, ..., vn is a basis for the vector space v.

By definition, a basis spans the entire vector space, which means every vector in v can be written as a linear combination of v1, v2, ..., vn.

Since the system is a basis, it must also be linearly independent, which implies that no vector in the system can be expressed as a linear combination of the other vectors.

Thus, the number of vectors in the system, n, is equal to the dimension of the vector space v, denoted as dim v.

If n = dim v, then the system of vectors is a basis:

Suppose n = dim v, where n is the number of vectors in the system and dim v is the dimension of the vector space v.

Since dim v is defined as the maximum number of linearly independent vectors that can form a basis for v, we know that any system of n linearly independent vectors in v will be a basis for v.

Therefore, the system of vectors v1, v2, ..., vn is a basis for the vector space v.

Combining both directions of the proof establishes that a linearly independent system of vectors v1, v2, ..., vn in a vector space v is a basis if and only if n = dim v.

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H does not satisfy all three properties required for a subspace, we can conclude that H is not a subspace of R2.

The set H is a subspace of R2, we need to check if it satisfies the three properties required for a subspace

1. The zero vector is in H.

2. H is closed under vector addition.

3. H is closed under scalar multiplication.

Now each property

1. The zero vector (0, 0) is in H since it lies on the unit circle.

2. To check closure under vector addition, suppose we have two vectors (x₁, y₁) and (x₂, y₂) in H. If we add them together, (x₁, y₁) + (x₂, y₂), the resulting vector will not necessarily lie on the unit circle. For example, if we add (1, 0) and (-1, 0), the result is (0, 0), which is not on the unit circle. Therefore, H is not closed under vector addition.

3. To check closure under scalar multiplication, suppose we have a scalar c and a vector (x, y) in H. If we multiply them, c × (x, y), the resulting vector will not necessarily lie on the unit circle. For example, if we multiply (1, 0) by 3, the result is (3, 0), which is not on the unit circle. Therefore, H is not closed under scalar multiplication.

Since H does not satisfy all three properties required for a subspace, we can conclude that H is not a subspace of R2.

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The solution to the system of equations is x = -67/3 and y = -17/9.

To solve the system of equations -x + 6y = 11 and 2x - 3y = 5, we can use the method of substitution or elimination. Let's use the elimination method to solve for x:

Multiply the first equation by 2 and the second equation by -1 to eliminate x:

-2(-x + 6y) = 2(11) --> 2x - 12y = 22

-1(2x - 3y) = -1(5) --> -2x + 3y = -5

Now, add the two equations together:

(2x - 12y) + (-2x + 3y) = 22 + (-5)

-9y = 17

Divide both sides of the equation by -9:

y = -17/9

Now, substitute the value of y back into one of the original equations. Let's use the first equation:

-x + 6(-17/9) = 11

-x - 34/3 = 11

Add 34/3 to both sides:

-x = 11 + 34/3

-x = 33/3 + 34/3

-x = 67/3

Multiply both sides by -1:

x = -67/3

Therefore, the solution to the system of equations is x = -67/3 and y = -17/9.

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The shear stress acting along the seam is 159.94 kPa.

We need to determine shear stress acting along the seam.

First, we are going to determine horizontal stress components at 0°. Then, using transformation formulas.

To find stress along the inclined seam.

We need to determine the cross-sectional area of the tube so we can calculate stress components.

A = π [tex](100/2)^{2}[/tex] -  [tex](100-4/2)^{2}[/tex]

= 615.8 [tex]mm^{2}[/tex]

= 6.16 * [tex]10^{-4} m^{2}[/tex]

Since the tube is only subjected to horizontal compressive force P at 0° there is only a normal stress component σ_x.

σ_x = P/A

​ σ_x = -200/6.16 * [tex]10^{-4}[/tex]

= -324806 Pa

Now we can apply the transformation formula for the shear stress component (9-2).

[tex]T_{x'}_{y'}[/tex] = - σ_x -  σ_y/2 sin 2θ + [tex]T_{xy}[/tex] cos 2θ

[tex]T_{x'}_{y'}[/tex] = -( -324806 -0/2) sin(2 * 40°) + 0 * cos(2 * 40°)

= 159936 Pa

≈ 159.94 kPa

Therefore [tex]T_{x'}_{y'}[/tex] = 159.94 kPa.

The Question was Incomplete, Find the full content below:

A paper tube is formed by rolling a cardboard strip in a spiral and then gluing the edges together as shown. Determine the shear stress acting along the seam, which is at 40° from the vertical when the tube is subjected to an axial compressive force of 200 N. The paper is 2mm thick and the tube has an outer diameter of 100 mm.

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The equation of the linear function is f(x) = 20x - 150.

To find the equation of the linear function, we need to find the slope and the y-intercept.

Using the given points, we can find the slope:

slope = (f(9) - f(6)) / (9 - 6) = (30 - (-30)) / 3 = 20

Now, to find the y-intercept, we can use one of the points. Let's use (6, -30):

y - y1 = m(x - x1)

y - (-30) = 20(x - 6)

y + 30 = 20x - 120

y = 20x - 150

Therefore, the equation of the linear function is f(x) = 20x - 150.

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To find out how many applications can be made with the 26-gram container, we need to divide the total amount of cream in the container by the average amount of cream used per application.

Total amount of cream (container) = 26 grams
Average amount of cream per application = 1 6/7 grams

First, let's convert the mixed fraction 1 6/7 to an improper fraction:
(1 * 7) + 6 = 13/7 grams

Now, divide the total amount of cream by the average amount of cream per application:

26 grams ÷ 13/7 grams

To divide by a fraction, you multiply by its reciprocal (the fraction flipped):

26 * 7/13

Now, cancel out the common factor (13):

(26/13) * (7/1)

2 * 7 = 14

So, you can make 14 applications with the 26-gram container.

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The dot plot that most likely models an unfair spinner is C. Seven marbles are above one. Two marbles are above two. Two marbles are above three. Two marbles are above four. Seven marbles are above five.

The only dot plot that is not likely to model a fair spinner is the third one. In this dot plot, 7 marbles land on the first sector, 2 marbles land on the second sector, 2 marbles land on the third sector, 2 marbles land on the fourth sector, and 7 marbles land on the fifth sector. This distribution is not likely to occur if the spinner is fair, as each sector should have an equal chance of landing face up.

The other three dot plots are more likely to model a fair spinner. In the first dot plot, each sector has 4 marbles land on it. In the second dot plot, each sector has 3 or 4 marbles land on it. In the fourth dot plot, each sector has 4 or 5 marbles land on it. These distributions are more likely to occur if the spinner is fair, as each sector has an equal chance of landing face up.

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To test for the significance of the coefficient on aggregate price index, we need to calculate the p-value.

The p-value is the probability of obtaining a result as extreme or more extreme than the one observed, assuming that the null hypothesis is true.

In this case, the null hypothesis would be that there is no relationship between the aggregate price index and the variable being studied. We can use statistical software or tables to determine the p-value.

Generally, if the p-value is less than 0.05, we can reject the null hypothesis and conclude that there is a significant relationship between the aggregate price index and the variable being studied. If the p-value is greater than 0.05, we cannot reject the null hypothesis.

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The length of side a is determined as 13.92 by applying sine rule of triangle.

The length of side a is calculated by applying the following formulas shown below;

Apply sine rule as follows;

a / sin (83) = 13 / sin (68)

Simplify the expression as follows;

multiply both sides of the equation by " sin (83)".

a = ( sin (83) / sin (68) ) x 13

a = 13.92

Thus, the value of side length a is determined as 13.92 by applying sine rule as shown above.

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The given initial-value problem is x(dy/dx)y = 2x + 1, y(1) = 9.

To solve this problem, we first rearrange the equation as (1/y) dy = (2/x + 1/x) dx. We can integrate both sides, which gives us ln|y| = 2ln|x| + ln|x| + b, where b is the constant of integration.

Simplifying this expression, we get ln|y| = 3ln|x| + b. Exponentiating both sides, we obtain |y| = eᵇ * x³. Since y(1) = 9, we substitute x = 1 and y = 9 into the equation, which gives us 9 = eᵇ * 1³, or b = ln 9. Therefore, the solution to the initial-value problem is y = ±9x³.

To solve this initial-value problem, we first rearranged the given equation to put it in a form that we can integrate. We then integrated both sides of the equation, introducing a constant of integration. By substituting the initial value of y, we were able to determine the value of the constant of integration and thus find the final solution to the initial-value problem.

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